Home Blog Page 5

Continuous functions of one variable

0

We discuss the important class of continuous functions of one variable. We teach you how to prove that a function is continuous. We also provide several examples and exercises with detailed answers to illustrate our course.

To fully understand the content of this page, some tools on the limit of functions are necessary.

Generalities on continuous functions

Throughout this paragraph,  $\mathscr{D}$ is a subset of $\mathbb{R}$ and $f:\mathscr{D}\to \mathbb{R}$ is a function. The set $\mathscr{D}$ is called the domain of $f$.

Definition: We say that $f$ is continuous at point $a\in D$ if $f$ admit a limit $f(a)$ at the point $a$. That is, for any $\varepsilon>0,$ there exists $\delta>0$, such that for any $x\in D,$ $|x-a|<\delta$ implies that $|f(x)-f(a)|<\varepsilon$.

We say that $f$ is continuous on $\mathscr{D}$ if $f$ is continuous on each point of $\mathscr{D}$.

If $f$ and $g$ are continuous functions, the sum, the product, the composition, and the quotient of these functions are also continuous.

Classical examples: All polynomial functions are continuous on $\mathscr{D}=\mathbb{R}.$ The square root function $x\mapsto \sqrt{x}$ is continuous on $\mathscr{D}=[0,+\infty)$. The trigonometric function $x\mapsto \sin(x)$ and $x\mapsto \cos(x)$ are continuous on $\mathbb{R}$. The logarithmic function $x\mapsto \ln(x)$ is continuous on $\mathscr{D}=(0,+\infty)$. The exponential function $x\mapsto e^x$ is continuous on $\mathscr{D}=\mathbb{R}$.

The other functions are usually the sum, product, composition, inverse, and quotient of the classical functions above.

Continuity using sequences: the function $f$ is continuous at $a\in \mathscr{D}$ if and only if for any sequence $(u_n)_n$ of $\mathscr{D}$ converging to $0,$ the image sequence $( f (u_n))_n$ converges to $f(a)$.

This result is very useful if we want to prove that the function $f$ is not continuous at a point $a$. To do this, it suffices to find two sequences both converging towards $a,$ while their images converge towards different points.

Theorem: If $f:[a,b]\to\mathbb{R}$ is continuous then $f$ is bounded on $[a,b]$, this is there exists a real number $M>0$ such that $|f(x)|\le M$ for any $x\in [a,b]$. Moreover, there exists $\alpha,\beta\in [a,b]$ such that \begin{align*} f(\alpha)=\sup_{x\in [a,b]}f(x),\qquad f(\beta)=\inf_{x\in [a,b]}f(x).\end{align*}

A selection of exercises on continuity of functions

Exercise: Determine if the following functions are continuous or not: \begin{align*} f(x)=\begin{cases}\displaystyle\frac{\sin(2x)-\sin(5x)}{\sin(3x)},& x\neq 0,\cr -1,& x=0.\end{cases}\qquad g(x)=\begin{cases} x^x,& x>0,\cr 2,&x=0.\end{cases} \end{align*}

Solution: 1) The function $f$ is continuous on $\mathbb{R}\backslash\{0\}$ as the sum and the quotient of continuous functions. Now let us verify the continuity of $f$ at $0$. For that purpose, we shall calculate the limit of $f$ as $0$ and check if it is equal to $f(0)$. Before doing so we recall the following standard result \begin{align*} \lim_{x\to 0,\;x\neq 0} \frac{\sin(x)}{x}=1. \end{align*} Now for any $x\neq 0,$ we have \begin{align*} f(x)&=\frac{x}{\sin(3x)}\times \left(\frac{\sin(2x)}{x}-\frac{\sin(5x)}{x}\right)\cr &= \frac{1}{3}\frac{3x}{\sin(3x)}\times \left(2\frac{\sin(2x)}{2x}-5\frac{\sin(5x)}{5x}\right). \end{align*} Hence \begin{align*} \lim_{x\to 0,\;x\neq 0}f(x)=\frac{1}{3}(2-5)=-1=f(0). \end{align*} It follows that $f$ is continuous on $\mathbb{R}$.

2) For any $x>0$ we can write \begin{align*} g(x)=e^{\ln(x^x)}=e^{x \ln(x)}. \end{align*} This shows that $g$ is continuous on $(0,\infty)$. On the other hand we now that \begin{align*} \lim_{x\to 0^+}x\ln(x)=0. \end{align*} This implies that \begin{align*} \lim_{x\to 0^+} g(x)=1\neq g(0). \end{align*} Then $g$ is not continuous at $0$.

Exercise: Consider the function \begin{align*} f(t)=\begin{cases}\frac{\ln(1+2t^2)}{\sin^2(t)} ,& t > 0,\cr 0,& x=2. \end{cases} \end{align*} Prove that $f$ is continuous on $[0,+\infty)$.

Solution: The function $f$ is continuous on $(0,+\infty)$, as the quotient of continuous functions. We now show what happens if $t$ is very closed to $0$. For $t\neq 0,$ we have \begin{align*} f(t)=2\frac{\ln(1+2t^2)}{2t^2}\left(\frac{t}{\sin(t)}\right)^2 \end{align*} As \begin{align*} \lim_{t\to 0,t\neq 0} \frac{\ln(1+2t^2)}{2t^2}=1, \end{align*} then \begin{align*} \lim_{t\to 0,t\neq 0}f(x)=2=f(0). \end{align*} This shows that $f$ is also continuous at $0$.

Exercise: Let $f$ be defined by \begin{align*} f(x)=\begin{cases} 1,& x\in \mathbb{Q},\cr 0,& x\notin \mathbb{Q}.\end{cases} \end{align*} Prove that $f$ is not continuous at any point of $\mathbb{R}$.

Solution: We recall that $g:A\subset \mathbb{R}\to \mathbb{R}$ is continuous at $x\in A$ if and only if for any sequence $(u_n)_n\subset A$ such that $u_n\to x$ as $n\to\infty,$ we have $g(u_n)\to g(x)$.

Using this result, to prove that $f$ is not continuous at $x\in \mathbb{R}$ it suffices to find two sequences $(u_n)_n$ and $(v_n)_n$ both converging to $x$ as $n\to \infty$ but the limits of $f(u_n)$ and $f(v_n)$ as $n\to\infty$ are different. In fact, let $x\in \mathbb{R}$. By density of $\mathbb{Q}$ in $\mathbb{R}$ there exists $(u_n)_n\subset \mathbb{Q}$ such that $u_n\to x$ as $n\to \infty$. By definition of $f$ we have $f(u_n)=1$, so the limit of $f(u_n)$ at infinity is $1$. On the other hand, by the density of $\mathbb{R}\backslash\mathbb{Q}$ in $\mathbb{R}$ there exists $(v_n)_n\subset \mathbb{R}\backslash\mathbb{Q}$ such that $v_n\to x$ as $n\to infty$. By definition of $f,$ we have $f(v_n)=0,$ so the limit of $f(v_n)$ at infinity is $0$. We then conclude that $f$ is not continuous at $x$.

How to find the eigenvalues of a matrix?

0

Eigenvalues of a matrix are a fundamental concept in linear algebra that provide valuable information about the behavior of matrices. They play a crucial role in various fields, including physics, engineering, data analysis, and computer science. In this article, we will explore the concept of eigenvalues, understand their significance, and delve into methods for finding the eigenvalues of a matrix.

We assume that you are familiarized with the matrix operations.

What are the eigenvalues of a matrix?

Eigenvalues of a matrix are scalar values associated with square matrices that reveal important properties of the matrices. For an $n \times n$ matrix $A$, an eigenvalue $\lambda$ is a value for which there exists a non-zero vector $x$, known as an eigenvector, such that $Ax = λx$. Eigenvalues and eigenvectors capture how matrices stretch, rotate, or scale vectors in different directions.

The characteristic equation of matrix

One way to find eigenvalues is by solving the characteristic equation. The characteristic equation for matrix A is given by $$P_A(\lambda):=\det(\lambda I-A):=|\lambda I-A| = 0, $$where I is the identity matrix. Solving this equation yields the eigenvalues of the matrix $A$.

Note that $P_A(\cdot)$ is a polynomial, called the characteristic polynomial of the matrix $A$.

The characteristic equation allows us to uncover the eigenvalues by finding the values of $\lambda$ that make the determinant of $\lambda I-A$ equal to zero.

The set of eigenvalues of a matrix $A$ is called the spectrum of $A$ and will be denoted by $\sigma(A)$. Some time to compute the eigenvalues of a matrix it suffices to solve simple quadratic equations.

The eigenspace associated with that eigenvalue

By definition, the eigenspace associated with an eigenvalue $\lambda$ of a matrix $A$ is the kernel of the matrix $\lambda I-A$, that is, $\ker(\lambda I-A)$. In the following example, we show how to determine the set of eigenvalues and the associated eigenspace.

Example

Determine the eigenvalues and their associated eigenspaces of the following matrix $$A=\begin{pmatrix} 2&-1&1\\ -1&2&-1\\ -1&1&0\end{pmatrix}.$$

For any $\lambda\in\mathbb{C}$ we have \begin{align*} \det( \lambda I_3-A)= \begin{vmatrix} \lambda-2&1&-1\\1&\lambda-2&1\\1&-1&\lambda\end{vmatrix}. \end{align*} We know that the determinant is unchanged if we replace any line of the matrix with a combination of the other lines “also if we replace any column with a linear combination of the others columns”. We denote by $L_i$ for $i=1,2,3,$ the lines of any matrix of order $3$. By replacing the line $L_1$ by $L_1+L_3$, we obtain \begin{align*} \det( \lambda I_3-A)&= \begin{vmatrix} \lambda-1&0&\lambda-1\\1&\lambda-2&1\\1&-1&\lambda\end{vmatrix}\cr &= (\lambda-1)\begin{vmatrix} 1&0&1\\1&\lambda-2&1\\1&-1&\lambda\end{vmatrix} \end{align*} In the lest determinant we replace the line $L_2$ by $L_2-L_3,$ we obtain \begin{align*} \det( \lambda I_3-A)&= (\lambda-1)\begin{vmatrix} 1&0&1\\0&\lambda-1&1-\lambda\\1&-1&\lambda\end{vmatrix}\cr &=(\lambda-1)^2\begin{vmatrix} 1&0&1\\0&1&-1\\1&-1&\lambda\end{vmatrix} \end{align*} We replace the line $L_3$ by $L_3+L_2-L_1$, we obtain \begin{align*} \det( \lambda I_3-A) &=(\lambda-1)^2\begin{vmatrix} 1&0&1\\0&1&-1\\0&0&\lambda-2\end{vmatrix}\cr &= (\lambda-1)^2(\lambda-2). \end{align*} Then the matrix $A$ process two eigenvalues “$\lambda_1=1$ is a double eigenvalue and $\lambda_2=2$ is a simple eigenvalue”.

Let us denote by $E_1$ the characteristic space associated with eigenvalue $\lambda_1=1$ and $E_2$ the characteristic space associated with eigenvalue $\lambda_2=2$. By definition we have \begin{align*} E_1=\ker(I_3-A),\quad E_2=\ker(2I_3-A). \end{align*} Then \begin{align*} X=\left(\begin{smallmatrix}x\\y\\z\end{smallmatrix}\right)\in E_1&\;\Longleftrightarrow\; A X=X\cr &\;\Longleftrightarrow\;\begin{cases} x-y+z=0\\ -x+y_z=0\\ -x+y-z=0\end{cases} \cr &\;\Longleftrightarrow\; x-y+z=0 \cr &\;\Longleftrightarrow\; X=\left(\begin{smallmatrix}y-z\\y\\z\end{smallmatrix}\right) \cr &\;\Longleftrightarrow\; X=\left(\begin{smallmatrix}y\\y\\0\end{smallmatrix}\right)+\left(\begin{smallmatrix}-z\\0\\z\end{smallmatrix}\right)\cr &\;\Longleftrightarrow\; X=y\left(\begin{smallmatrix}1\\1\\0\end{smallmatrix}\right)+z\left(\begin{smallmatrix}-1\\0\\1\end{smallmatrix}\right)\cr &;\Longleftrightarrow\; X\in {\rm span}\left\{\left(\begin{smallmatrix}1\\1\\0\end{smallmatrix}\right),\left(\begin{smallmatrix}-1\\0\\1\end{smallmatrix}\right)\right\}. \end{align*} This shows that \begin{align*} E_1={\rm span}\left\{\left(\begin{smallmatrix}1\\1\\0\end{smallmatrix}\right),\left(\begin{smallmatrix}-1\\0\\1\end{smallmatrix}\right)\right\}. \end{align*} Now we calculate $E_2$. Let $X=\left(\begin{smallmatrix}x\\y\\z\end{smallmatrix}\right)$. Then \begin{align*} X\in E^2&\;\Longleftrightarrow\; AX=2X\cr &\;\Longleftrightarrow\;\begin{cases} -y+z=0\\ -x-z=0\\-x+y-2z=0\end{cases} \cr &\;\Longleftrightarrow\; y=z,\quad x=-z \cr &\;\Longleftrightarrow\; X=\left(\begin{smallmatrix}-z\\z\\z\end{smallmatrix}\right)\cr &\;\Longleftrightarrow\; X=z\left(\begin{smallmatrix}-1\\1\\1\end{smallmatrix}\right)\cr &\;\Longleftrightarrow\; X\in {\rm span}\left\{\left(\begin{smallmatrix}-1\\1\\1\end{smallmatrix}\right)\right\}. \end{align*} Hence \begin{align*} E_2={\rm span}\left\{\left(\begin{smallmatrix}-1\\1\\1\end{smallmatrix}\right)\right\}. \end{align*}

Algebraic Multiplicity and Geometric Multiplicity

Eigenvalues can have both algebraic multiplicity and geometric multiplicity. The algebraic multiplicity represents the number of times an eigenvalue appears as a root of the characteristic equation, while the geometric multiplicity reflects the dimension of the eigenspace associated with that eigenvalue. Understanding the relationship between these multiplicities provides insights into the behavior of the matrix.

In the previous example, we have seen that the matrix $A$ has two eigenvalues $\lambda_1=1$ and $\lambda_2=2$. Clearly, the algebraic multiplicity of $\lambda_1$ is equal to 1, while the algebraic multiplicity of $\lambda_2$ is 2. We have also proved that the eigenspace $\ker(I-A)$ is spanned by a non-null vector, that $\dim(\ker(I-A))=1$. Thus the geometric multiplicity of the eigenvalue $\lambda_1=1$ is equal to $1$. On the other hand, we have also proved that the eigenspace $\ker(I-A)$ is spanned by two linearly independent vectors. Therefore, $\dim(\ker(2I-A))=2$. Thus the geometric multiplicity of the eigenvalue $2$ is equal to $2$.

Computing Eigenvalues

Several methods can be employed to find eigenvalues efficiently. These include power iteration, which iteratively calculates dominant eigenvalues, and the QR algorithm, which provides a systematic approach to finding all eigenvalues. Additionally, specialized algorithms such as the Jacobi method and the Arnoldi iteration are used for large and sparse matrices.

Applications of Eigenvalues of a Matrix

Eigenvalues have broad applications in different fields. In physics, they help understand the behavior of physical systems, such as quantum mechanics and oscillatory systems. In data analysis and machine learning, eigenvalues are crucial for dimensionality reduction techniques like principal component analysis (PCA). Eigenvalues are also essential in solving differential equations and studying dynamic systems.

Conclusion Eigenvalues of a matrix

Eigenvalues provide valuable insights into the behavior of matrices and have extensive applications across various fields. By understanding the concept of eigenvalues and employing methods for their computation, we gain a deeper understanding of the underlying structure and properties of matrices. The ability to find and analyze eigenvalues empowers us to tackle complex problems, make accurate predictions, and unlock the potential of linear algebra in diverse domains.

The optimal step gradient method

0

In this article, we propose an exercise that describes the optimal method of the step gradient for coercive functions on spaces of finite dimension. These specific problems are used in optimization theory and some applications in finance. We mention also we can use the convex functions to study the optimality of some models.

A problem with the optimal method of the step gradient

Problem: Let $f:\mathbb{R}^n\to \mathbb{R}$ be a function of class $C^1$. Suppose that there exists a real number $\alpha>0$ such that \begin{align*} \forall (u,v)\in \mathbb{R}^n\times \mathbb{R}^n,\qquad \langle \nabla f(v)-\nabla f(u),v-u\rangle \;\ge\alpha |v-u|^2. \end{align*}

  • Prove that \begin{align*} \forall (u,v)\in \mathbb{R}^n\times \mathbb{R}^n,\qquad f(v)\ge f(u)+ \langle\nabla f(u),v-u\rangle+\frac{\alpha}{2} \|v-u\|^2. \end{align*}
  • Deduce that $f$ admits a global minimum on $\mathbb{R}^n$, reached in a unique point that will be denoted by $a$.
  • Let $u\in \mathbb{R}^n\backslash\{a\}$. Consider the function \begin{align*}\varphi_u:\mathbb{R}\to \mathbb{R},\quad t\mapsto f(u+t\nabla f(u)). \end{align*} Prove that $\varphi_u$ admits a global minimum on $\mathbb{R},$ reached in a unique point. This allows us to define a sequence $(u_k)_{k\ge 0}$ in $\mathbb{R}^n$ such that: $u_0\in \mathbb{R}^n$ ; if $k_k=a,$ we select $u_{k+1}=a,$ if not, let $t_k$ be the unique real number such that $\varphi_{u_k}(t_k)=\min(\varphi_{u_k})$. We then set \begin{align*} u_{k+1}=u_k+t_k \nabla f(u_k). \end{align*}
  • Verify that for any $k\in\mathbb{N},$ $\nabla f(u_k)\bot \nabla f(u_{k+1})$.
  • Prove that the series $(u_{k+1}-u_k)_{k\ge 0}$; $(\nabla f(u_{k+1})-\nabla f(u_k))_{k\ge 0}$ and $(\nabla f(u_k))_{k\ge 0}$ the suites tend towards $0$. Deduce that $u_k\to a$ as $k\to\infty$.

Hyperplane linear subspace

0

We will delve into the concept of a hyperplane linear subspace, explore its properties, and examine its applications in different domains.

Hyperplanes are essential concepts in linear algebra that provide a framework for understanding linear subspaces in multidimensional spaces. They play a crucial role in various fields, including geometry, machine learning, and optimization.

What is the hyperplane linear space?

Let us start with the following definition.

Definition of the hyperplane

A hyperplane is a subspace in a higher-dimensional space that has one less dimension than the space itself.

In a two-dimensional space, a hyperplane is simply a line. In three-dimensional space, a hyperplane is a two-dimensional plane. Generally, in an n-dimensional space, a hyperplane is an (n-1)-dimensional subspace.

If $E$ is a vector space of finite dimension equal to $n\ge 2$, then for any hyperplane $H$ of $E$, we have $\dim(H)=n-1$.

Remark

Let $E$ be a vector space such that $\dim(E)=n\ge 2$ and $f:E\to \mathbb{R}$ be a linear form on $E$. According to the rank theorem $\dim(\ker(f))=n-1$. Thus $\ker(f)$ is a hyperplane of $E$. In addition if we denote by $\langle\cdot,\cdot\rangle$ the canonique product space in $E$, the there exists a unique $a\in E$ such that $f(x)=\langle x,a\rangle$.

If $H$ is an hyperplane of $E$ then there exists a unique $a\in E$ such that $$ H=\{x\in E: \langle x,a\rangle=0\}.$$

The equation of a hyperplane 

Let $E$ be a vector space of dimension $n\ge 2$ on a field $\mathbb{R}$. Thus $E$ can be viewed as the space $\mathbb{R}^n$. We then define a product space of $E$ by: for vectors $x=(x_1,\cdots,x_n)$ and $y=(y_1,\cdots,y_n)$ in $E$, with $x_i$ and $y_i$ are real numbers, then $$ \langle x,y\rangle=x_1y_1+\cdots+x_ny_n.$$ According to the above remark, if $H$ is an hyperplane of $E$, then there exists a unique $(a_1,\cdots,a_n)\in E\setminus 0$ such that $ x\in H$ if and only if $ \langle x,a\rangle=0$. thus the hyperplanes are the solutions of the equations of the form $$ a_1x_1+\cdots+a_n x_n=0.$$ Observe that the hyperplane is normal to the vector $(a_1,\cdots,a_n)$.

Exercises on the hyperplane linear subspace

Here we propose some exercises with detailed solutions to illustrate the concept of a hyperplane linear space.

Exercise 1

Let $V$ be a vector space and $n\in\mathbb{N}$, let $H$ be a hyperplane of $V$, and let $v\in V$ be a vector. Under what condition the subspaces $H$ and ${\rm span}(v):=\{\lambda v:\lambda\in\mathbb{C}\}$ are supplementally in $V$.

We shall discuss two cases: First case: if $v\in H$. Then for any $\lambda\in\mathbb{C},$ $\lambda v\in H$. Thus ${\rm span}(v)\subset H,$ so that $H$ and ${\rm span}(v)$ are not supplementally. Second case: if $v\notin H$. First of all, we have $\dim({\rm span}(v))=1$. As $H$ is a hyperplane of $V,$ it follows that $\dim(H)=n-1$. Hence \begin{align*} \dim(H)+\dim({\rm span}(v))=n=\dim(V). \end{align*} Now let $x\in H\cap {\rm span}(v)$. Then $x\in H$ and there exists $\lambda\in\mathbb{C}$ such that $x=\lambda v$. We necessarily have $\lambda=0,$ because if not, then $v=\lambda^{-1} x\in H,$ absurd. Hence $x=0_V$, and then $H\cap {\rm span}(v)={0_V}$. This shows that $H+{\rm span}(v)=V$ is a direct sum.

Exercise on hyperplanes 2

Let $\psi: \mathbb{R}^n\to \mathbb{C}$ be a nonnull linear forme and $\Phi$ be an endomorphism of $\mathbb{R}^n$. Prove that the kernel of $\psi$ is stable by $\Phi$, i.e. $\Phi(\ker(\psi))\subset \ker(\psi)$, if and only if there exists a real $\lambda\in\mathbb{R}$ such that $\psi\circ\Phi=\lambda\psi$.

Assume that there exists $\lambda\in\mathbb{R}$ such that $\psi\circ\Phi=\lambda\psi$. Let $x\in \ker(\psi)$ and prove that $\Phi(x)\in \ker(\psi)$. In fact, we have $\psi(x)=0$. On the other hand, \begin{align*} \psi(\Phi(x))=\lambda \psi(x)=0. \end{align*} This implies that $\Phi(\ker(\psi))\subset \ker(\psi)$. Conversely, assume that $\ker(\psi)$ is stable by $\Phi$. Observe that if $x\in \ker(\psi)$ then $\psi(\Phi(x))=0,$ so that \begin{align*} \psi(\Phi(x))=\lambda \psi(x),\quad \forall x\in \ker(\psi),\;\forall \lambda\in\mathbb{R}. \end{align*} It suffice then to look for real $\lambda$ and a supplementally space $K$ of $\ker(\psi)$ such that $\psi\circ\Phi=\lambda\psi$ on $K$. According to rank theorem we have $\ker(\psi)$ is a hyperplane, so $\dim(\ker(\psi))=n-1$. Thus any supplementally space $K$ of $\ker(\psi)$ will satisfies $\dim(K)=1$. Take $a\in \mathbb{R}^n$ such that $\psi(a)\neq 0$, a such $a$ exists because $\psi$ is a non null forme, so that $a\notin \ker{\psi}$. This implies that ${\rm span}(a)\cap \ker(\psi)=\{0\}$. As $\dim({\rm span}(a))=1$ and $\dim(\ker{\psi})+\dim({\rm span}(a))=n$. Then \begin{align*} \ker(\psi)\oplus {\rm span}(a)=\mathbb{R}^n. \end{align*} Then it suffices to look for real $\lambda$ such that $\psi\circ\Phi=\lambda\psi$ on ${\rm span}(a)$. In particular $\psi(\Phi(a))=\lambda \psi(a)$. We choose \begin{align*} \lambda=\frac{\psi(\Phi(a))}{\psi(a)}. \end{align*}

Properties of Hyperplanes

Hyperplanes possess several important properties that make them significant in linear algebra. Firstly, they divide the space into two regions, with points on one side satisfying the inequality $$a_1x_1 + a_2x_2 + \cdots + a_nx_n \le b,$$ and points on the other side satisfying the inequality $$a_1x_1 + a_2x_2 + \cdots + a_nx_n \ge b.$$ Additionally, the normal vector of a hyperplane is orthogonal to any vector lying within the hyperplane. This orthogonality property has applications in various mathematical and computational algorithms.

Applications of Hyperplanes

Hyperplanes find diverse applications in multiple fields. In geometry, hyperplanes help define convex sets and facilitate the study of convex hulls and polytopes. In machine learning and pattern recognition, hyperplanes are instrumental in supporting vector machines (SVM) and separating data points with different class labels. Optimization algorithms often utilize hyperplanes to define constraints and boundaries for solving optimization problems. Hyperplanes also play a role in computational geometry, computer graphics, and data visualization.

Hyperplane Intersection and Dimensionality Reduction

The intersection of hyperplanes can reveal interesting properties and relationships. When hyperplanes intersect in an n-dimensional space, the resulting subspace has a dimension less than (n-1). This intersection property can be leveraged in dimensionality reduction techniques such as Principal Component Analysis (PCA) to capture the most significant features and reduce the complexity of high-dimensional datasets.

Conclusion on hyperplane linear subspace

Hyperplanes form a fundamental concept in linear algebra, serving as linear subspaces that define boundaries and separate spaces in multidimensional settings. Understanding hyperplanes and their properties enables us to comprehend geometric relationships, perform classification tasks, and solve optimization problems. The versatility of hyperplanes is evident in their applications across diverse fields, ranging from machine learning to computational geometry. By exploring the concept of hyperplanes, mathematicians, data scientists, and researchers can unlock the potential of linear subspaces, paving the way for innovative solutions and advancements in various disciplines.

Elementary probability exercises

0

We offer elementary probability exercises at the level of high school. We show the student how to use the formula of compound probabilities, how to use conditional probability, and the Bayes theorem.

What is a random experiment?

To study phenomena, it is important to carry out experiments. In fact, depending on the situation, certain phenomena are governed by deterministic laws, and the associated experiences have only one possible result. For other phenomena, however, governed by laws, the experiments lead to several results, forming all the possible results. Thus, two similar experiments can lead to very different results. We cannot predict the exact outcome of this experiment, so it is a random experiment.

We mention that there is a deep relationship between probability and certain physical phenomena such as heat equations. Likewise where the solution is associated with a probability law called Gaussian distribution. So that it is the advanced probability that is not the subject of this article.

In this post, we only discuss elementary probability properties.

Elementary probability exercises

Exercise: An urn contains 3 white balls and 5 red balls. We make four drawn without discount. Determine the probability of firing four red balls.

Solution: We denote by $R_i$ the event ” the i-th drawn ball is red”. Our objective is to calculate the probability of the event $A=R_1\cap R_2\cap R_3\cap R_4$.

It is extremely difficult to calculate each $\mathbb{P}(R_i)$, but the conditional probabilities $\mathbb{P}(R_i|R_1\cdots R_{i-1})$ are simple to calculate. In fact, the realization of the event $R_1\cdots R_{i-1}$ shows that the urn contains 3 white balls and $5-i+1$ red balls. Then we have a probability equal to $\frac{6-i}{9-i}$ of sorting a red ball. Hence, the formula of compound probabilities gives \begin{align*} \mathbb{P}(A)&=\mathbb{P}(R_1)\mathbb{P}(R_2|R_1)\mathbb{P}(R_3|R_1\cap R_2)\mathbb{P}(R_4|R_1\cap R_2\cap R_3)\cr &= \frac{5}{8}\times \frac{4}{7}\times \frac{3}{6}\times \frac{2}{5}\cr &= \frac{1}{14}. \end{align*}

Exercise: Steve and James train in archery. Steve reaches the target 9 times out of 10, James only 6 times out of 10. James plays two out of three. To reach the target what is the probability?

Solution: We denote by $S$ the event ” Steve plays”, $J$ the event “James plays”, and $T$ the event “the target is reached”. As $J=\overline{S}$, $(S,J)$ is a complete system of events. The formula of total probabilities implies \begin{align*} \mathbb{P}(T)&=\mathbb{P}(T|S)\mathbb{P}(S)+\mathbb{P}(T|J)\mathbb{P}(J)\cr &= \frac{9}{10}\times \frac{1}{3}+\frac{6}{10}\times \frac{2}{3}\cr &= \frac{7}{10}. \end{align*}

Riemann Integral Exercises

0

Riemann integral exercises with detailed answers are offered on this page. These kinds of integrals of bounded function in bounded intervals, almost everywhere continuous. Thus, all continuous or monotone functions are Riemann integrable.

Lower integral, upper integral, and proper integral 

In this section, we give a concise definition of Riemann integrals. To this end, we assume that a real function is defined and bounded on the interval $[a,b]$, with $a,b\in\mathbb{R}$ are such that $a<b$.

For a subdivision $\sigma:=\{x_0,\cdots.x_n\}$ of the interval $[a,b]$, i.e. $a=x_0<x_1<\cdots<x_n=b$, we define the lower sum of $f$ on $[a,b]$, by $$ L_n(f,\sigma)=\sum_{i=1}^n m_i(f,\sigma)(x_i-x_{i-1}),$$ where $m_i(f,\sigma)$ is the infinimum of $f$ on the interval $[x_{i-1},x_i]$.

By definition, the lower integral of f over the interval $[a,b]$ is $$ \underline{\int^b_a}f(x)ds=\sup_{\sigma} L_n(f,\sigma)=\lim_{n\to\infty}U_n(f,\sigma).$$ Similarly, we define the upper sum of $f$ over the interval $[a,b]$ by $$ U_n(f,\sigma)=\sum_{i=1}^n M_i(f,\sigma)(x_i-x_{i-1}),$$ where $M_i(f,\sigma)$ is the supremum of $f$ on the interval $[x_{i-1},x_i]$.

The upper integral of f over the interval $[a,b]$ is $$ \overline{\int^b_a}f(x)ds=\inf_{\sigma} U_n(f,\sigma)=\lim_{n\to\infty}U_n(f,\sigma).$$ Definition: A bounded function $f:[a,b]\to\mathbb{R}$ is said to be Riemann integrable on $[a,b]$ if its lower and upper integrals coincide. In this case,  the integral of $f$ over $[a,b]$ is given by $$\int^b_a f(x)dx=\overline{\int^b_a}f(x)ds=\underline{\int^b_a}f(x)ds.$$ There exists another equivalent condition for Riemann integrals. Let’s give a summary of it. We take $c_\in [x_{i-1},x_i]$ for $i=1,\cdots,n$ and de define the Riemann sum by $$S_n(f,\sigma):=\sum_{i=1}^n f(c_i)(x_{i}-x_{i-1}).$$ Observe that $L_n(f,\sigma)\le S_n(f,\sigma)\le U_n(f,\sigma)$. Thus if $f$ is Riemann integrable then $S_n(f,\sigma)$ has a limit as $n\to\infty,$ and in this case we have $$ \int^b_a f(x)dx=\lim_{n\to\infty} \sum_{i=1}^n f(c_i)(x_{i}-x_{i-1}).$$

A particular cas: Let $f:[0,1]\to \mathbb{R}$ be Riemann integral and take the uniform subdivision $\sigma:=\{\frac{i}{n}: i=0,\cdots,n\}$, and $c_i=\frac{i}{n}$, Then $$ \int^1_0 f(x)dx=\lim_{n\to\infty}\frac{1}{n}\sum_{i=1}^n f\left(\frac{i}{n}\right).$$

Riemann integral exercises

We start our selection of Riemann integral exercises by providing an example of a non-integrable bounded function in the sense of Riemann

Exercise:  Prove that the following function $$g(x)=\cos^2(x)1_{\mathbb{Q}}(x),\quad x\in [0,\pi/2],$$ is not Riemann integrable.

Exercise: Let $f:[0,1]\to \mathbb{R}$ be a continuous function. For any $n\in\mathbb{N},$ we set \begin{align*} u_n=n\int^1_0 t^n f(t)dt. \end{align*}

  1. Let $\varepsilon>0$. Prove that there exists $a\in (0,1)$ such that $|f(t)|\le \frac{\varepsilon}{2}$ whenever $t\in [a,1]$.
  2. We denote $M=\sup{|f(t)|:t\in [0,1]}$, a such $M\in \mathbb{R}^+$ exists because $f$ is continuous on the compact $[0,1]$. Prove that \begin{align*} |u_n|\le M a^n+\frac{\varepsilon}{2},\qquad \forall n\in\mathbb{N}. \end{align*}
  3. Deduce that $\lim_{n\to+\infty}u_n=0$.
  4. Deduce that, in general, we have \begin{align*} \lim_{n\to+\infty}u_n=f(1). \end{align*}

Solution: The first 3 question: By definition of left continuity at point $1,$ for all $\varepsilon>0,$ there exists $0 < \delta < 1$ such that $t\in [1-\delta,1]$ we have $|f(t)-f(1)|=|f(t)| \le \frac{\varepsilon}{2}$. It suffices to put $a=1-\delta\in (0,1)$.

We can write \begin{align*} |u_n|&\le n \int^1_0 t^n |f(t)|dt\cr &\le n \int^a_0 t^n |f(t)|dt+n \int^1_a t^n |f(t)|dt\cr & \le n M \int^a_0 t^n dt+ n \frac{\varepsilon}{2} \int^1_a t^n dt\cr & \le n M \left[\frac{t^{n+1}}{n+1}\right]^a_0+ n \frac{\varepsilon}{2} \left[\frac{t^{n+1}}{n+1}\right]^1_a \cr &\le M \frac{n}{n+1} a^{n+1}+\frac{\varepsilon}{2} \frac{n}{n+1} (1-a^{n+1}). \end{align*} As $a\in (0,1),$ then $a^{n+1}\le a^n$ and $0 < 1-a^{n+1} < 1$. On the other hand, $\frac{n}{n+1}le 1$. This implies that \begin{align*} |u_n|\le M a^n+\frac{\varepsilon}{2},\qquad \forall n\in\mathbb{N}. \end{align*}

Since $a\in (0,1)$ then the geometric sequence $(a^n)_n$ satisfies $a^n\to 0$ as $n\to+\infty$. Then there exists $N\in \mathbb{N}$ such that for any $n\ge N,$ we have $0 < a^n < \frac{\varepsilon}{2M}$. Hence for $n\ge N,$ we have \begin{align*} |u_n|\le M \frac{\varepsilon}{2M}+\frac{\varepsilon}{2}=\varepsilon. \end{align*} Thus $u_n\to 0$ as $n\to +\infty$.

4) In this question we work with a general continuous function $f:[0,1]\to \mathbb{R}$. We set $g(t)=f(t)-f(1)$. The function $g$ is continuous on $[0,1]$ and satisfies $g(1)=0$. According to the first question we have \begin{align*} \lim_{n\to+\infty} n\int^1_0 t^n g(t)dt=0. \end{align*} In addition, observe that \begin{align*} u_n&=n\int^1_0 t^n g(t)dt+n f(1) \int^1_0 t^n dt \cr &= n\int^1_0 t^n g(t)dt+\frac{n}{n+1}f(1) \end{align*} Hence $u_n\to f(1)$ as $n\to +\infty$.

Exercise: Let $f:\mathbb{R}\to \mathbb{R}$ be a continuous function and let $\omega\in\mathbb{R}^\ast$. Define the following function \begin{align*} \varphi(x)=\frac{1}{\omega}\int^x_0 \sin\left(\omega (x-t)\right)f(t),dt,\qquad x\in\mathbb{R}. \end{align*}

  • Prove that there exist two functions $\Phi,\Psi:\mathbb{R}\to \mathbb{R}$ such that for any $x\in\mathbb{R},$ \begin{align*} \varphi(x)=\frac{\sin(\omega x)}{\omega} \Phi(x)- \frac{\cos(\omega x)}{\omega}\Psi(x). \end{align*}
  • Show that $\varphi$ is twice differentiable on $\mathbb{R}$ and that \begin{align*} \varphi”+\omega^2 \varphi=f. \end{align*}

Solution: 1) Let $x\in\mathbb{R}$ and $t\in\mathbb{R}$. We know from trigonometric formulas that \begin{align*} \sin(\omega x-\omega t)=\sin(\omega x)\cos(\omega t)-\cos(\omega x)\sin(\omega t). \end{align*} Then \begin{align*} \varphi(x)&= \frac{\sin(\omega x)}{\omega} \int^x_0 \cos(\omega t)f(t)dt-\frac{\cos(\omega x)}{\omega} \int^x_0 \sin(omega t)f(t)dt. \end{align*} Thus it suffices to select \begin{align*} \Phi(x)= \int^x_0 \cos(\omega t)f(t)dt\quad\text{and} \quad \Psi(x)=\int^x_0 \sin(\omega t)f(t)dt. \end{align*}

2) The functions $\Phi$ and $\Psi$ are primitives of continuous functions. Thus $\Phi$ and $\Psi$ are $C^1$ functions on $\mathbb{R}$ and that \begin{align*} \Phi'(x)=\cos(\omega x)f(x),\qquad \Psi(x)=\sin(\omega x)f(x) \end{align*} for any $x\in \mathbb{R}$. This proves that the function $varphi$ is differential on $\mathbb{R}$ as the product and sum of differentiable functions. Moreover, \begin{align*} \omega\varphi'(x)&= \omega \cos(\omega x) \Phi(x)+\sin(\omega x) \cos(\omega x)f(x)\cr & \hspace{2cm}+\omega\sin(\omega x)\Psi(x)-\cos(\omega x)\sin(\omega x)f(x)\cr &= \omega \cos(\omega x) \Phi(x)+\omega\sin(\omega x)\Psi(x). \end{align*} As $\omega\neq 0,$ then \begin{align*} \varphi'(x)=\cos(\omega x) \Phi(x)+\sin(\omega x)\Psi(x) \end{align*} This last formula show that $\varphi’$ is differentiable as product and sum of differentiable functions. Hence $\varphi$ is twice differentiable and \begin{align*} \varphi”(x)&=-\omega\sin(\omega x) \Phi(x)+\cos^2(\omega x)f(x)\cr & \hspace{2.5cm}+\omega\cos(\omega x)\Psi(x)+\sin^2(\omega x)f(x)\cr &= -\omega (\sin(\omega x) \Phi(x)-\cos(\omega x)\Psi(x))+f(x)\cr &= -\omega^2 \varphi(x)+f(x) \end{align*} for any $x\in\mathbb{R}$. This proves that \begin{align*} \varphi”+\omega^2 \varphi=f. \end{align*}

Translation operator

0

We discuss the properties of the translation operator defined in a Lebesgue space. The latter appears in the study of differential equations, mainly equations with negative memories called also delay equations.

The semigroup property of the translation operator 

Let $p\ge 1$ be a real number and $E=L^p([0,+\infty)$ be the space of all measurable functions such that \begin{align*} \int^{+\infty}_0 |f(s)|^pds <\infty.\end{align*} If we select \begin{align*} \|f\|_p:=\left( \int^{+\infty}_0 |f(s)|^pds <\infty \right)^{\frac{1}{p}}\end{align*}\then $(E,\|\cdot\|)$ is a Banach space of infinite dimension, the Lebesgue space. On the other hand, let us define \begin{align*} (T_t f)(s)=f(t+s),\qquad \forall s,t\ge 0,\; f\in E.\end{align*}

Fact 1: for any $t\ge 0,$ $T_t$ is a linear continuous operator on $E,$ i.e. $T_t\in \mathscr{L}(E)$. Moreover, $\|T_t\|_{\mathscr{L}(E)}\le 1$ for any $t$.

Proof: Let $t\ge 0$ and $f\in E$. The function $s\mapsto (T_t f)(s)=f(t+s)$ is measurable because as composition of two measurables functions, $f$ and the function $s\mapsto t+s$. Moreover, we estimate \begin{align*} \int^{+\infty}_0 |(T_tf)(s)|^pds&= \int^{+\infty}_0 |f(t+s)|^pds \cr & =\int_t^{+\infty}|f(\sigma)|^pd\sigma\cr & \le \int^{+\infty}_0 |f(\sigma)|^pd\sigma=\|f\|_p^p.\end{align*} This shows that $T_f f\in E$, and \begin{align*} \|T_t f\|_p\le \|f\|_p,\qquad \forall t\ge 0,\;f\in E.\end{align*}

Definition: The operator $T_t$ is called the translation operator or the shift operator.

Fact 2: $T_0=Id_E$ and for any $t,s\ge 0,$ $T_{t+s}=T_tT_s$, here $T_tT_s:=T_t\circ T_s$.

Proof: For $f\in E$ and $t,s\ge 0,$ we have \begin{align*} (T_{t+s}f)(s)=f(s+(t+\sigma))=(T_sf)(t+\sigma)=(T_t T_sf)(\sigma).\end{align*}Definition: The family of operators $(T_t)_{t\ge 0}$ satisfying the fact 1 and the fact 2 is called the semigroup of operators or the shift semigroup.

The strong continuity of the shift semigroup 

In this section, we show that the translation operator $T_t$ satisfies a nice topological property.

Fract 3: The family $(T_t)_{t\ge 0}$ is strongly continuous at zero. That is for any $f\in E,$ we have $$\lim_{t\to 0}\|T_t f-f\|_p=0.$$

Proof:  For $f\in E,$ we select $u(t):=T_tf$. According to Fact 2, it suffices to prove that $\|u(t)-u(0)\|_p\to 0$ as $t\to 0$. Before solving this question, we first need to recall some facts. First, a function $f$ has a support compact if it vanishes outside a compact. The support of $f$ is the compact set \begin{align*}{\rm supp}(f):=\overline{\{t: f(t)\neq 0\}}.\end{align*} Second, we recall that the set of continuous functions with compact support, $\mathscr{C}_c([0,+\infty))$, is dense in $L^p([0,+\infty))$.

As $f\in E,$ then by density there exists a sequence $(f_n)_n\subset \mathcal{C}_c([0,+\infty)) $ such that $\|f_n-f\|_p\to 0$ as $n\to+\infty$. For any $\varepsilon,$ there exists $m\in \mathbb{N}$ such that $\|f_m-f\|\le \frac{\varepsilon}{3}$. On the other hand, we have\begin{align*}\|u(t)-u(0)\|&=\|T_t f-f\|_p\cr & \le \|T_t f-T_t f_m\|_p + \|T_t f_m-f_m\|_p+\|f_m-f\|_p\cr & \le 2 \|f_m-f\|_p + \|T_t f_m-f_m\|_p \cr &\le \frac{2}{3}\varepsilon+ \|T_t f_m-f_m\|_p .\end{align*}

Now it suffices to show that $\|T_t f_m-f_m\|_p\to 0$ as $t\to 0$. In fact, put $K={\rm supp}(f)$. Then, \begin{align*} \|T_t f_m-f_m\|_p ^p&=\int_K |f_m(t+s)-f(s)|^p ds\cr &\le {\rm meas}(K) \left(\sup_{s\in K}|f(t+s)-f(s)|\right )^p .\end{align*} Here $ {\rm meas}(K) $ is the Lebesgue measure of the set $K$. As the function $f$ is uniformly continuous on the compact set, then we have the following uniform convergence \begin{align*}\lim_{t\to 0} \sup_{s\in K}|f(t+s)-f(s)| =0.\end{align*} Then there exists $\delta>0$ such that for any $t\in (0,\delta)$, we have \begin{align*} \|T_t f_m-f_m\|_p \le \frac{\varepsilon}{3}.\end{align*}Thus, for any $\varepsilon,$ there exists $\delta>0$ such that for any $t,$ \begin{align*} t\in (0,\delta)\Longrightarrow \|u(t)-u(0)\| < \varepsilon.\end{align*}

Properties of Real Numbers

0

Real numbers are the backbone of mathematics, spanning a wide range of values along the number line. They’re like the superheroes of numbers, including familiar integers, fractions, and decimals, both rational and irrational. But did you know that the properties of real numbers are actually used in real life? That’s right, these numbers represent quantities, distances, and magnitudes in the real world, making them an essential concept in mathematics. From simple addition and subtraction to complex calculus, these versatile numbers are the key to unlocking the mysteries of our universe’s quantitative aspects. So, if you want to be a math whiz, start by mastering the power of real numbers!

Integers

Hey there! If you’re a math whiz, you probably already know about counting numbers. But just in case you need a refresher, let’s talk about integers. Integers are whole numbers – that means they can be positive, negative, or even zero! We use the symbol $\mathbb{Z}$ to represent the set of all integers.

Now, let’s move on to whole numbers. These are non-negative integers, which means they can’t be negative but they can be zero. We use the symbol $\mathbb{N}$ to represent the set of all whole numbers.

Last but not least, let’s talk about natural numbers. These are positive integers, which means they can’t be zero but they can be any positive number. We use the symbol $\mathbb{N}^\ast$ to represent the set of all natural numbers.

So there you have it – a quick rundown of counting numbers, integers, whole numbers, and natural numbers. Happy counting!

What is a real number?

Real numbers are the backbone of mathematics, representing positions along the number line and encompassing a vast range of values. They come in two main categories: rational and irrational numbers. Rational numbers are those that can be expressed as a fraction of two integers, including integers, fractions, and repeating or terminating decimals: $$ \mathbb{Q}=\left\{\frac{p}{q}: p,q\in\mathbb{Z},\;q\neq 0\right\}.$$ On the other hand, irrational numbers are those that cannot be expressed as a simple fraction and have non-repeating, non-terminating decimal expansions. We use the symbol $\mathbb{Q}^c$ to represent the set of all irrational numbers. With this notation, the set of real numbers is $$ \mathbb{R}=\mathbb{Q}\cup \mathbb{Q}^c.$$

Finally, non-negative real numbers include all real numbers greater than or equal to zero, including positive real numbers, zero, and all non-positive real numbers.

Real numbers are used extensively in mathematics to represent quantities, measurements, and values in various scientific disciplines. They are a fundamental concept in mathematical analysis and are essential for understanding the world around us. So, whether you’re a math whiz or just starting to explore the world of numbers, real numbers are a crucial concept to grasp.

Properties of real numbers

In this paragraph, we list some important properties of real numbers.

Closure property :

Did you know that real numbers have a superpower? They possess the closure property, which means that when you add or multiply two real numbers, the result will always be another real number. It’s like they have a secret code that keeps them together, making them closed under addition and multiplication. Pretty cool, right? If $a$ and $b$ are real numbers, the $a+b$ and $ab$ are also real numbers,<

Commutative Property :

The commutative property is a fundamental concept in mathematics that pertains to the addition and multiplication of real numbers. Specifically, it denotes that the order in which numbers are arranged does not impact the outcome of these operations. In the case of addition, the commutative property is expressed as $a + b = b + a$, while for multiplication, it is represented as $a b = b a$.

Associative Property:

The Associative Property pertains to the operations of addition and multiplication of real numbers, whereby the grouping of numbers does not alter the outcome. Specifically, in the case of addition, the expression $(a + b) + c$ is equivalent to $a + (b + c)$, while in the case of multiplication, $(a b)c$ is equivalent to $a (b c)$.

Identity Property:

The identity element for addition is denoted as $0$, as it satisfies the condition that the sum of any real number $a$ and $0$ is equal to $a$. Similarly, the identity element for multiplication is represented by the number $1$, as it satisfies the condition that the product of any real number $a$ and $1$ is equal to $a$. This property is essential in various mathematical applications and serves as a basis for further mathematical operations.

Inverse Property:

For any real number $a$, there exists an additive inverse, denoted as $(-a)$, such that the sum of $a$ and its additive inverse is equal to zero, i.e., $a + (-a) = 0$. Furthermore, for any nonzero real number $a$, there exists a multiplicative inverse, denoted as $\frac{1}{a}$, such that the product of $a$ and its reciprocal is equal to one, i.e., $a \times \frac{1}{a} = 1$. This property is of great significance in various mathematical fields, including algebra, calculus, and number theory.

Distributive Property

The distributive property is a fundamental mathematical principle that asserts that multiplication is distributive over addition. Specifically, it states that the product of a given number, denoted by “$a$,” and the sum of two other numbers, denoted by “$b$” and “$c$,” is equivalent to the sum of the products of “$a$” and each of the two numbers, “$b$” and “$c$.” This can be expressed mathematically as $a \times (b + c) = (a \times b) + (a \times c)$.

Transitive Property

The Transitive Property is a fundamental principle in mathematics that states that if two quantities, $a$ and $b$, are equal, and $b$ and $c$ are also equal, then $a$ and $c$ must be equal as well. This property plays a crucial role in establishing equalities and inequalities involving real numbers, and is widely employed in various mathematical disciplines.

Ordering Property

The concept of ordering property pertains to the arrangement of real numbers in ascending or descending order. Specifically, if a real number $a$ is less than another real number $b$, it is denoted as $a < b$. Conversely, if $a$ is greater than $b$, it is represented as $a > b$. The concept of ordering property pertains to the arrangement of real numbers in ascending or descending order. Specifically, if a real number a is less than another real number b, it is denoted as a < b. Conversely, if a is greater than b, it is represented as a > b.

Density Property

The Density Property is a fundamental characteristic of the real number line, which asserts that an infinite number of real numbers exist between any two given real numbers. This property underscores the continuity of the real number line. Specifically, for any two real numbers $a$ and $b$ such that $a < b$, there exists at least one rational number $r$ that satisfies the inequality $a < r < b$.

The aforementioned properties serve as the fundamental basis for resolving equations, inequalities, and diverse mathematical operations that entail real numbers. Their significance in the fields of algebra and calculus is paramount, as they facilitate the comprehension and manipulation of real-world quantities and data by mathematicians and scientists.

What is the Absolute Value of a Real Number?

The absolute value of a real number, represented by the symbol |x| and pronounced as “the absolute value of x,” is a metric that quantifies the distance between the number and zero on the number line. This metric provides a measure of the magnitude of a number, irrespective of its sign. The absolute value is always a non-negative real number or zero. The definition of the absolute value is as follows:

  • For a positive real number x, the absolute value is equivalent to the number itself: $|x| = x$.
  • For a negative real number x, the absolute value is the additive inverse of that number: $|x| = -x$.
  • For zero, the absolute value remains zero: $|0| = 0$.

We recall the following triangular inequality: for all real numbers $x$ and $y$, $$ |x+y|\le |x|+|y|.$$

Remarkable identities

The utilization of the following noteworthy identities is deemed advantageous:

    \begin{align*} &(a+b)^2=a^2+2ab+b^2\cr & (a-b)^2=a^2-2ab+b^2\cr & (a+b)(a-b)=a^a-b^2\cr & (a+b)^n=\sum_{k=0}^n \binom{n}{k} a^kb^{n-k},\end{align*} where $$ \binom{n}{k}=\frac{n!}{k!(n-k)!},\quad n!=1\times 2\times 3\times\cdots\times n.$$

These identities play a vital role in mathematics. The first two aid in expanding quadratic expressions, while the third facilitates the factorization of the difference between two squares. The fourth identity, known as the binomial theorem, allows for the expansion of binomial expressions raised to positive integer powers. In the fourth identity, the summation notation signifies that it involves a sum of terms. Each term results from the multiplication of a binomial coefficient by the product of the appropriate powers of the two terms in the binomial expression.

What is an integer part of a real number?

The “floor” of a real number, commonly known as its integer part, is defined as the largest integer that is less than or equal to the given real number. This operation effectively truncates the decimal portion of the number, leaving only the integer component. The notation used to denote the integer part of a real number x is commonly represented as $[x]$. The integer part of a real number satisfies the following important property

For all real number, $$ [x]\le x < [x]+1.$$
For instance, if we consider the real number $5.7$, its integer part is $[5.7] = 5$, since the largest integer that is less than or equal to $5.7$ is $5$. Similarly, for the real number $-3.2$, its integer part is $[-3.2]= -4$, as the largest integer that is less than or equal to $-3.2$ is $-4$. In the case of $8.0$, since there is no decimal component, its integer part is $[8.0]= 8$. Finally, for the irrational number $\pi,$ approximately $3.14159$, its integer part is $[\pi] = 3$, as the largest integer that is less than or equal to $\pi$ is $3$.

Frequently Asked Questions on Real Numbers

What Are Real Numbers? +
Real numbers are a set of numbers that includes all rational and irrational numbers. They can be represented on the number line and encompass a wide range of values, including integers, fractions, decimals, and more.
What Are Rational Numbers? +
Rational numbers are real numbers that can be expressed as fractions, where the numerator and denominator are integers, and the denominator is not zero. Examples include 1/2, -3/4, and 7.
What Are Irrational Numbers? +
Irrational numbers are real numbers that cannot be expressed as fractions and have non-repeating, non-terminating decimal expansions. Examples include $\sqrt{2} and $\pi$.
What Is the Absolute Value of a Real Number? +
The absolute value of a real number, denoted as $|x|$, represents the distance of that number from zero on the number line. It is always a non-negative value. For example, $|5|=5$ and $|-5|=5$.
What Are the Properties of Real Numbers? +
Real numbers are like a treasure trove of mathematical properties! They’re closed, meaning that when you add or multiply them, you always end up with another real number. Plus, they’re super friendly – addition and multiplication are commutative and associative, so you can rearrange and group them however you like. And let’s not forget about their trusty sidekicks, the identity elements – 0 for addition and 1 for multiplication. But wait, there’s more! Real numbers also have inverse elements, which means that for every number, there’s a special friend that, when added or multiplied together, gives you the identity element. And finally, the cherry on top – the distributive property, which lets you distribute multiplication over addition “or vice versa” like a boss. Who knew numbers could be so cool?

Linear transformations and matrices

0

Linear transformations and matrices form the backbone of numerous mathematical and scientific disciplines, ranging from computer graphics and machine learning to physics and engineering. The interplay between these concepts is not only intellectually stimulating but also crucial in solving real-world problems efficiently.

In this article, we will explore the deep connection between linear transformations and matrices, shedding light on their fundamental principles and practical applications.

Generalities on linear transformations

To begin, let’s define what linear transformations and matrices are. A linear transformation is a function that preserves vector addition and scalar multiplication. In simpler terms

Definition of a linear transformation

Let $E$ and $F$ be vector spaces on a field $\mathbb{K}=\mathbb{R}$ or $\mathbb{C}$. An applicaion $f: E\to F$ is called a linear transformation if for any $u,v\in E$ and $\lambda\in\mathbb{K},$ \begin{align*} f(u+\lambda v)=f(u)+\lambda f(v).\end{align*}

The set of all linear transformations from $E$ to $F$ is denoted by $\mathcal{L}(E,F)$ and it is a vector space.

The space $\mathcal{L}(E):=\mathcal{L}(E,E)$ is called the space of endomorphism.

Examples of linear spplications

$\bullet$ Let $E$ be a vector space and $\lambda\in \mathbb{R}$. Define the application $f:E\to E$ by $f(x)=\lambda x$. This $f$ is a linear transformation.

$\bullet$ Let $T:\mathbb{R}^3\to \mathbb{R}^2$ the application defined by $$ T\begin{pmatrix} x\\y\\z\end{pmatrix}=\begin{pmatrix}2x-y+z\\ y-z\end{pmatrix}.$$ Then $T$ is a linear application.

$\bullet$ Let $R[X]$ be the space of all polynomial with real coefficients and let $g: \mathbb{R}[X]\to \mathbb{R}$ be defined by $$ g(P)=P(0)+P'(0),\qquad P\in \mathbb{R}[X].$$ Then $g$ is a linear application.

The first property of linear transformations $f$ is $f(0)=0$. In fact, $f(0)=f(0+0)=f(0)+f(0)=2f(0)$, which implies that $f(0)=0$.

The range of a linear transformation $f: E\to F$ denoted as ${\rm Im}(f)$ is defined by $${\rm Im}(f)=\{f(u):u\in E\}.$$

The kernel of a linear map $f: E\to F$ is denoted as $\ker(f)$ and defined by $$\ker(f):=\{u\in E:f(u)=0\}.$$ It is a subspace of $E$.

We recall that an application $f: E\to F$ is injective if for $x,y\in E$ with $f(x)=f(y)$ implies that $x=y$. Now if in addition $f$ is linear, then the injectivity of $f$ is equivalent to $\ker(f)=\{0\}$.

The rank theorem:

If the spaces $E$ and $F$ have a finite dimension and if the map $f: E\to F$ is linear, then \begin{align*} \dim(E)=\dim(\ker(f))+\dim({\rm Im}(f)).\end{align*} By the way the number $\dim({\rm Im}(f))$ is called the rank of $f$ and will be denoted by ${\rm rank}(f)$

Here we gather some classical examples of linear transformations with detailed solutions.

Linear transformations examples

Example on the space of continuous functions

Let $E=\mathscr{C}(\mathbb{R},\mathbb{R})$ be the vector space on $\mathbb{R}$ of all continuous functions from $\mathbb{R}$ to $\mathbb{R}$. If $f$ and $g$ are real functions then the usual product of $f$ by $g$ is defined by $(fg)(x)=f(x)g(x)$ for any $x\in\mathbb{R}$. We denote by $id:\mathbb{R}\to \mathbb{R}$ the identity function, $id(x)=x$ for all $x\in \mathbb{R}$. Let us now define the application $\Phi:E\to E$ by \begin{align*} \Phi(f)(x)=xf(x),\qquad \forall f\in E,\quad x\in \mathbb{R}. \end{align*}

Prove that $\Phi$ is a linear transformation and determine the kernel of $\Phi$. Does $\Phi$ surjective? Determine the image of $\Phi$.

First of all, we need to prove that the map $\Phi$ is well defined. In fact, if $f\in E$, the the function $x\mapsto xf(x)$ is continuous as product of continuous functions. This implies that $\Phi(f)\in E$, so $\Phi$ define a function from $E$ to $E$. Let us now prove that $\Phi$ is a linear transformation on $E$. In fact, let $f,g\in E$ and $\lambda\in \mathbb{R}$. For any $x\in \mathbb{R},$ we have \begin{align*} \Phi(f+\lambda g)(x)&= x(f(x)+\lambda g(x))\cr &= xf(x)+\lambda (xg(x))\cr &=\Phi(f)(x)+\lambda \Phi(g)(x). \end{align*} Hence \begin{align*} \Phi(f+\lambda g)=\Phi(f)+\lambda \Phi(g). \end{align*} This means that $\Phi$ is linear.

Next, let us compute the the kernel of $\Phi$. For $f\in \ker(\Phi)$, we have $xf(x)=0$ for any $x\in \mathbb{R}$. This equality always hold for $x=0$. Now assume that $x\neq 0$, this implies that $f(x)=0$. But as $f$ is continuous at zero, by taking the limit of $f$ at zero we obtain $f(0)=0$. this means that $f$ is, in fact, identically null on $\mathbb{R}$. Thus $f=0$, and then $\ker(\Phi)=\{0\},$ which means that the application $\Phi$ is injective.

The map $\Phi$ is not surjective. In fact, let define the function $g(x)=1$ for all $x\in\mathbb{R}$, so that $g\in E$ as a constant function. Now if we assume that $\Phi$ is surjective then we can find $f\in E$ such that $\Phi(f)=1$. This means that for any $x\in \mathbb{R}$ we have $xf(x)=1$. But if we take $x=0$, we find $0=1$. This is absurd, and hence $\Phi$ is not surjective.

Finally, let us determine ${\rm Im}(\Phi)$. Let $g\in {\rm Im}(\Phi)$. Then there is $f\in E$ such that $g(x)=xf(x)$ for all $x\in \mathbb{R}$. In particular, $g(0)=0$. On the other have \begin{align*} \forall x\in \mathbb{R}\setminus\{0\},\quad \frac{g(x)-g(0)}{x-0}=f(x). \end{align*} As $f$ is continuous at $0,$ we deduce that $g$ is differentiable in $0$ and $g'(0)=f(0)$. This implies that \begin{align*} {\rm Im}(\Phi)\subset \{g\in E: g(0)=0\;\text{and}\; f\;\text{is differentiable at}\; 0\}. \end{align*} Conversely, let $g\in E$ such that $g(0)=0$ and $g$ is differentiable at $0$. Then the function \begin{align*} f(x)=\begin{cases}\frac{g(x)}{x},& x\neq 0,\cr g'(0),& x=0,\end{cases} \end{align*} is an element of $E$ and $g=\Phi(f)\in {\rm Im}(\Phi)$. Hence \begin{align*} {\rm Im}(\Phi)= \{g\in E: g(0)=0\;\text{and}\; f\;\text{is differentiable at}\; 0\}. \end{align*}

The following example uses linear transformations to solve a problem on supplementary spaces.

Example on linear transformations and supplementary spaces:

Let $\mathbb{K}$ be a field and consider the $\mathbb{K}$-vector space $E=\mathbb{K}^n$ with $n\ge 0$. Let $a=(1,1,\cdots,1)\in E$ and denote \begin{align*} A&:=\mathbb{K}a={(\lambda,\lambda,\cdots,\lambda):\lambda\in \mathbb{K}}\cr B&:=\{x=(x_1,x_2,\cdots,x_n)\in E: x_1+x_2+\cdots+x_n=0\}. \end{align*} Let the linear transformation \begin{align*} f:E\to \mathbb{K},\quad f(x_1,x_2,\cdots,x_n)=\sum_{i=1}^n x_i. \end{align*}

  • Determine $\ker(f)$. Deduce that $B$ is a subspace.
  • Let $x\in E$. Prove that there exists a unique scalar $\lambda\in\mathbb{K}$ such that $f(x)=\lambda f(a)$.
  • Justify that $A\cap B=\{0\}$.
  • Deduce that the sum $A+B$ is direct, that is, $E=A\oplus B$.

$\bullet$ The kernel of $f$ is, by definition, \begin{align*} \ker(f)&=\left\{(x_1,x_2,\cdots,x_n)\in E: f(x_1,x_2,\cdots,x_n)=0\right\}\cr &  =\left\{(x_1,x_2,\cdots,x_n)\in E: x_1+x_2+\cdots+x_n=0\right\}\cr &= B. \end{align*} As $B$ coincides with the kernel of a linear transformation, it is a subpace.

$\bullet$ Let’s prove the uniqueness of $\lambda$. Assume that there exist $(\lambda,\mu)\in\mathbb{K}^2$ such that $f(x)=\lambda f(a)=\mu f(a)$. As $f(a)=n,$ we then have $(\lambda-\mu) n=0$, so that $\lambda=\mu$. Now let us prove its existence. Let $x=(x_1,x_2,\cdots,x_n)\in X$. It suffice to prove that there exist $\lambda\in\mathbb{K}$ such that $x-\lambda a\in \ker(f)$, so that \begin{align*} \sum_{i=1}^n x_i-\lambda n=0. \end{align*} It suffice to take \begin{align*} \lambda=\frac{x_1+x_2+\cdots+x_n}{n}. \end{align*}

$\bullet$ Let $x\in A\cap B$ Then we have there exist $\lambda\in \mathbb{K}$ such that \begin{align*} x=\lambda a\quad\text{and}\quad x_1+x_2+\cdots+x_n=0. \end{align*} As $f$ is linear then $f(x)=\lambda f(a)$. By the previous question we have \begin{align*} \lambda=\frac{x_1+x_2+\cdots+x_n}{n}=\frac{0}{n}=0. \end{align*} Hence $x=\lambda a=0\times a=0$, so that $A\cap B=\{0\}$.

$\bullet$ Because of the previous question it suffices to show that $E=A+B$. According to the question $2$, there exists $\lambda\in\mathbb{K}$ such that $f(x)=\lambda f(a)$, so that $f(x-\lambda a)=0$. Then $x-\lambda a\in \ker(f)=B$. On the other hand, $\lambda a \in A$. Thus $x=\lambda a+(x-\lambda a)\in A+B$. This ends the proof.

The link between Linear Transformations and Matrices

Let $U$ and $V$ be two finite-dimensional vectors spaces with $\dim(U)=n$ and $\dim(V)=m$ and consider a linear transformation $T\in \mathcal{L}(U,V)$. Consider $\mathscr{B}_U=\{u_1,\cdots,u_n\}$ as the basis of $U$, and $\mathscr{B}_V=\{v_1,\cdots,v_m\}$ as the basis of $V$.

For any $i=1,2,\cdots,n$, we have $T(u_i)\in V$. Then we can represent the vector $T(u_i)$ in the basis $\mathscr{B}_V$ as $$ T(u_i)=\sum_{k=1}^m a_{ki}v_k,$$ wher $a_{ik}$ are scalars.

Then we associate to $T$ a matrix denoted as $$ {\rm Mat}(T):={\rm Mat}(T,\mathscr{B}_U\to \mathscr{B}_V)$$ and define by $$ {\rm Mat}(T)=\begin{pmatrix} T(u_1)&T(u_2)&\cdots & T(u_n)\end{pmatrix}$$ That is $$ {\rm Mat}(T)=\begin{pmatrix} a_{11}& a_{12}&\cdots & a_{1n}\\ a_{21}& a_{22}&\cdots & a_{2n}\\ \vdots&\vdots&\ddots&\vdots\\ a_{m1}& a_{m2}&\cdots & a_{mn}\end{pmatrix}.$$ We denote by ${\rm Mat}(m\times n)$ the vector space of all matrices with $m$ lines and $n$ columns. In addition if $A\in {\rm Mat}(m\times n)$, we dnote $A=(a_{ij})_{1\le i\le m,1\le j\le n}$.

Product of matrices

Consider three vector spaces $U,V$ and $W$ with bases $\mathscr{B}_U,\mathscr{B}_V$, and $\mathscr{B}_W$, respectively. We assume that $\dim(U)=n$, $\dim(V)=m$ and $\dim(W)=k$. Let $T\in\mathcal{L}(U,V)$ and $S\in\mathcal{L}(V,W)$. Then $$ S\circ T\in \mathcal{L}(U,W).$$ Let us denote $$ A={\rm Mat}(T,\mathscr{B}_U\to \mathscr{B}_V)\in {\rm Mat}(m\times n)$$ and $$ B={\rm Mat}(S,\mathscr{B}_V\to \mathscr{B}_W)\in {\rm Mat}(k\times m).$$ Then $$ C:=BA={\rm Mat}(S\circ T,\mathscr{B}_U\to \mathscr{B}_W)\in {\rm Mat}(k\times n)$$ The entries of the matrix $C$ are given by

$$ c_{ij}=\sum_{p=1}^m a_{ip}b_{pj}.$$

We know that the basis of any finite-dimensional vector space is not unique. Thus on the same vector space, one can find serval bases. Now have the following natural question:

Question

Consider $\mathscr{B}_U$ as the basis of $U$, and $\mathscr{B}_V$ as the basis of $V$. Let now $\mathscr{B}’_U$ and $\mathscr{B}’_V$ be two other bases of $U$ and $V$, respectively. Do we have $${\rm Mat}(T,\mathscr{B}_U\to \mathscr{B}_V)={\rm Mat}(T,\mathscr{B}’_U\to \mathscr{B}’_V)?.$$

The answer is no. Let us see this on simple example. Consider the linear application $T:\mathbb{R}^2\to \mathbb{R}^2$ defined by $$ T\begin{pmatrix}x\\ y\end{pmatrix}=\begin{pmatrix}x+y\\ x-y\end{pmatrix}.$$ Consider the vectors $$ e_1=\begin{pmatrix}1\\ 0\end{pmatrix},\quad e_2=\begin{pmatrix}0\\ 1\end{pmatrix},\quad v_1=\begin{pmatrix}1\\ 1\end{pmatrix}.$$ It is not diffuclt to prove that $\mathscr{B}=\{e_1,e_2\}$ and $\mathscr{B}’=\{e_1,v_2\}$ are bases of $\mathbb{R}^2$. Observe that $$ T(e_1)=e_1+e_2,\quad T(e_2)=e_1-e_2.$$ Then $$ {\rm Mat}(T,\mathscr{B}\to \mathscr{B})=\begin{pmatrix} 1&1\\1&-1\end{pmatrix}.$$ On the other hand, $$ T(e_1)=\begin{pmatrix}1\\1\end{pmatrix}=v_2,\qquad T(v_2)= \begin{pmatrix}2\\0\end{pmatrix}=2 e_1.$$ Then $$ {\rm Mat}(T,\mathscr{B}’\to \mathscr{B}’)=\begin{pmatrix} 0&2\\1&0\end{pmatrix}.$$

Lyapunov stability for nonlinear systems

0

We discuss the Lyapunov stability for nonlinear systems. Indeed, having a reference solution, a stationary solution, one wonders if the ODE solution is closed to this reference when the time is very large.

Throughout this post, we suppose that $F:\Omega\to \mathbb{R}^d$ is a continuously differentiable function, $C^1(\Omega)$, where $\Omega$ is an open set of $\mathbb{R}^d$. Moreover, we consider the Cauchy problem \begin{align*}\tag{CP}\dot{u}(t)=F(u(t)),\quad t\ge 0,\quad u(0)=x_0.\end{align*}

According to Peano’s, see also Cauchy-Lipschitz theorem, the Cauchy problem $(CP)$ admits a maximal solution.

The definition of exponentially stable equilibrium

We say that $x^\ast\in \Omega$ is an equilibrium, or critical, point of the application $F$ if it satisfies $F(x^\ast)=0$. Of course, a such point is not necessarily unique.

If $F$ is a linear transformation then one of the equilibrium points of $F$ is $0$.

If we define the constant function $u^\ast(t)=x^\ast$ for any $t$. Then we have $\dot{u}^\ast(t)=0=F(x^\ast)=F(u^\ast(t))$. This shows that $u^\ast$ is a solution of the differential equation $\dot{u}^\ast(t)=F(u^\ast(t))$. This solution is called the stationary solution, reference solution, or target solution.

Definition: We say that an equilibrium point $x^\ast\in \Omega$ is exponentially stable if There exist constants $omega, M>0$ and there exists $\delta>0$ such that if the initial state $u(0)=x_0$ satisfies $\|x_0-x^\ast\|\le \delta,$ then the maximal solution $u$ is golable, this means that $u(t)$ defined for any $t\in [0,+\infty)$, and we satisfies the following estimates \begin{align*}\|u(t)-x^\ast\|\le Me^{-\omega t}\|x_0\|,\qquad \forall t\ge 0.\end{align*}

This means that if the initial state $x_0$ is closed to the critical point, the solution is closed exponentially to this equilibrium.

Lyapunov’s stability theorem on nonlinear systems

We notice that without losing generalities, we can assume that the point of equilibrium of $F$ is $x^\ast$. This is because, if $x^\ast$ is an equilibrium point, then by replacing $u(t)$ by $v(t)=u(t)-x^\ast$ and $F(x)$ by $G(x)=F(x+x^\ast),$ then we have $G(0)=0$.

Theorem: Let $F$ be a $C^1(\Omega)$ class function such that $F(0)=0$, $0$ is a critical point of $F$. We denote by $\sigma(F'(0))$ the spectrum of the matrix $F'(0)$, the set of eigenvalues. We assume that the spectral bound \begin{align*} s(F'(0)):=\sup{{\rm Re}\lambda: \lambda\in \sigma(F'(0))} < 0. \end{align*} Then $0$ is exponentially stable.

Matrix Operations

0

Matrix operations are mathematical operations performed on matrices, which are rectangular arrays of numbers. Matrices are widely used in various fields, including mathematics, physics, computer science, and engineering.

What is a matrix?

Formally, a matrix is ​​just an array in which we put entries. This table can be a square or a rectangle.

Mathematically, a matrix is a linear transformation $A$ from $\mathbb{R}^n$ to $\mathbb{R}^p$. If $p=n$ then $A$ is called a square matrix, and if $p\neq q,$ $A$ is called a rectangular matrix. Using the base of the vector space $\mathbb{R}^p,$ the matrix $A$ takes the following form \begin{align*}A=\begin{pmatrix} a_{11}&a_{21}&\cdots&a_{n1}\cr a_{12}&a_{22}& \cdots&a_{n2}\cr \vdots&\vdots&\ddots&\vdots\\ a_{1p}&a_{2p}&\cdots&a_{np}\end{pmatrix}.\end{align*}For simplicity, we write $A=(a_{ij})_{\underset{1\le j\le p}{1\le i\le n}}$. We say that $A$ is a $np$-matrix. On the other hand, $n$ is the number of columns, while $p$ is the number of lines in the matrix.

Here are some common matrix operations:

In this section, we gather the essential matrix operations that are used every day!!!

Matrix Addition

Two matrices of the same size can be added together by adding their corresponding elements. For example, if A and B are both m×n matrices, the sum C = A + B is obtained by adding each element of A to the corresponding element of B.

If $$ A=\begin{pmatrix} 1&1&1\\1&1&1\\1&1&1\end{pmatrix},\quad B=\begin{pmatrix} 2&2&2\\2&2&2\\2&2&2\end{pmatrix},$$ then $$C=A+B=\begin{pmatrix} 3&3&3\\3&3&3\\3&3&3\end{pmatrix}.$$

Matrix Subtraction

Similar to addition, two matrices of the same size can be subtracted by subtracting their corresponding elements. If A and B are both m×n matrices, the difference C = A – B is obtained by subtracting each element of B from the corresponding element of A.

Matrix Scalar Multiplication

A matrix can be multiplied by a scalar, which is a single number. Each element of the matrix is multiplied by the scalar value. If A is an m×n matrix and c is a scalar, the product C = cA is obtained by multiplying each element of A by c.

If $$ A=\begin{pmatrix} 1&2\\ 3&4\end{pmatrix}$$ then $$ 3A=\begin{pmatrix} 3&6\\ 9&12\end{pmatrix}.$$

In what follows the field $\mathbb{K}$ will be the set of real numbers $\mathbb{R}$ or $\mathbb{C}$ the set of complex numbers.  We denote by $\mathscr{M}_{np}(\mathbb{K})$ the set of all matrices $A=(a_{ij})_{\underset{1\le j\le p}{1\le i\le n}}$. If $n=p,$ we set $\mathscr{M}_{nn}(\mathbb{K}):=\mathscr{M}_{n}(\mathbb{K}),$ the set of square matrices of order $n$.

Let $A$ and $B$ be two matrices in $\mathscr{M}_{np}(\mathbb{K})$ with coefficients $a_{ij}$ and $b_{ij},$ respectively. Let $\lambda \in \mathbb{K}$. We define the following matrix operations \begin{align*} A+B&=(a_{ij}+b_{ij})_{\underset{1\le j\le p}{1\le i\le n}}\cr \lambda\cdot A &=(\lambda a_{ij})_{\underset{1\le j\le p}{1\le i\le n}}.\end{align*} Then $(\mathscr{M}_{np}(\mathbb{K}),+,\cdot)$ is a vector space on $\mathbb{K}$.

Let’s consider the elementary matrices $E^{ij}$  for $i=1,\cdots,n$ and $j=1,\cdots,p$ defined in the following way: all the coefficients are zero except the coefficient which corresponds to index $i,j$. We recall that the set $$\{E^{ij}:i=1,\cdots,n,\;j=1,\cdots,p\}$$ is basis of the matrix space $\mathscr{M}_{np}(\mathbb{K})$. Then we have the following important result:

The dimension of the matrix vector space

$$ \dim(\mathscr{M}_{np}(\mathbb{K}))=np.$$

Matrix Multiplication is one of the complicated matrix operations

Matrix multiplication is a more complex operation that combines the rows of the first matrix with the columns of the second matrix to produce a new matrix. If A is an m×n matrix and B is an n×p matrix, the product C = AB is an m×p matrix. The (i, j)-th element of C is computed by taking the dot product of the i-th row of A and the j-th column of B.

Let matrices $A=(a_{ij})\in \mathscr{M}_{np}(\mathbb{K})$ and $B=(b_{ij})\in \mathscr{M}_{qn}(\mathbb{K})$. Then $C=AB=(c_{ij})\in \mathscr{M}_{pq}(\mathbb{K})$, where the entries $c_{ij}$ is given by \begin{align*}c_{ij}=\sum_{k=1}^n a_{ik}b_{kj}.\end{align*}

The power of a matrix $A$ is a matrix $$A^n=\underset{(n\;\text{times})}{A\cdot A\cdots A},\quad n\in\mathbb{N}.$$

When multiplying matrices you have to be very careful because in general $AB\neq BA$. In fact, takes, for example, matrices \begin{align*} A=\begin{pmatrix} 1&2\\1&0\end{pmatrix},\quad B=\begin{pmatrix}0&3\\ 1&2\end{pmatrix}.\end{align*} Then $$ AB=\begin{pmatrix}2&7\\ 0&3\end{pmatrix},\quad BA=\begin{pmatrix}3&0\\3&2\end{pmatrix}.$$ Then we cannot immediately use the binomial expansion formula since this formula is used in a commutative ring. Here the matrix space is not commutative. But we can still use the binomial expansion if the matrices $A$ and $B$ satisfy $AB=BA$.

Matrix operations: Matrix Transposition

The transposition of a matrix is obtained by interchanging its rows with columns. If $A$ is an m×n matrix, the transpose $A^t$ is an n×m matrix. The (i, j)-th element of $A^t$ is equal to the (j, i)-th element of $A$.

If $$ A=\begin{pmatrix} 1&5&4\\3&2&7\\0&1&8\end{pmatrix},$$ then $$ A^t=\begin{pmatrix} 1&3&0\\5&1&0\\4&7&1\end{pmatrix}.$$

 Matrix of a linear map

Let $\varphi: E\to F$ be a linear map, where $E$ and $F$ are finite-dimensional spaces with dimensions $n$ and $p$, respectively. That is $\varphi(x+\lambda y)=\varphi(x)+\lambda \varphi(y)$ for any $x,y\in E$, and $\lambda\in\mathbb{K}$. Let $(e_1,\cdots,e_n)$ and $(f_1,f_2,\cdots,f_p),$ the basies of $E$ and $F,$ respectively. The $np$-matrix $A$ associated with the linear map $\varphi$ is given by \begin{align*}A=\begin{pmatrix}\varphi(e_1)&\varphi(e_2)&\cdots&\varphi(e_n)\end{pmatrix}\end{align*}where for each $i=1,2,\cdots,n,$ $\varphi(e_i)$ is a colum vector calculted in the basis $(f_1,f_2,\cdots,f_p),$.

Example: Let $\mathbb{R}_2[X]$ the vector space of a polynomial of degree less or equal to $2$. We recall that the dimension of this space is $3$. Consider the application \begin{align*} \Phi:\mathbb{R}_2[X]\to \mathbb{R}^3,\qquad \Phi(P)=(P(-1),P(0),P(1)). \end{align*} Determine the matrix associated with $\Phi$. Prove that $\Phi$ is an isomorphism.

$\bullet$ Let us fisrt prove that $\Phi$ is a linear transformation. In fact, let $P,Q\in \mathbb{R}_2[X]$ and $\lambda\in \mathbb{R}$. We have \begin{align*} \Phi(P+\lambda Q)&=(P(-1)+\lambda Q(-1),P(0)+\lambda Q(0),P(1)+\lambda Q(1))\cr &= (P(-1),P(0),P(1))+(\lambda Q(-1),\lambda Q(0),\lambda Q(1)) \cr &= (P(-1),P(0),P(1))+\lambda (Q(-1), Q(0), Q(1))\cr &= \Phi(P)+\lambda \Phi(Q). \end{align*}

$\bullet$ Let $B=(1,X,X^2)$ the canonical basis of $\mathbb{R}_2[X]$ and $B’=(e_1,e_2,e_3)$ the canonical basis of $\mathbb{R}^3$. Now determine the matrix $A$ which represents the linear map in these bases. To compute the matrix associated with $\Phi$ in the basis $B$ et $B’$, we first give the coordinates of the vectors $\Phi(1),\Phi(X)$ and $\Phi(X^2)$ in the basis $B’$. We recall that $e_1=(1,0,0),\;e_2=(0,1,0),$ and $e_3=(0,0,1)$. Then \begin{align*} \Phi(1)&=(1,1,1)=e_1+e_2+e_3\cr \Phi(X)&=(-1,0,1)=-e_1+e_3\cr \Phi(X^2)&=(1,0,1)=e_1+e_3. \end{align*} Then \begin{align*} A&=\begin{pmatrix} \Phi(1)&\Phi(X)&\Phi(X^2)\end{pmatrix}\cr &= \begin{pmatrix} 1&-1&1\\1&0&0\\1&1&1\end{pmatrix}. \end{align*}

$\bullet$ Finally, we prove that $\Phi$ is an isomorphism. In fact, we know that \begin{align*}\dim\left(\mathbb{R}_2[X]\right)=3=\dim\left(\mathbb{R}^3\right).\end{align*} Then to prove that $\Phi$ is an isomorphism it suffices to prove that $\Phi$ is injective, which is equivalent to proving that the kernel $\ker(\Phi)={0}$. If $P\in \ker(\Phi)$ then $P(0)=P(-1)=P(1)=0$. This means that $P$ has three distinct roots, which is impossible because the degree of $P$ is less or equal $2$. Hence $P=0$ is “the null polynomial”, so $\ker(\Phi)={0}$. The application $\Phi$ is then injective then it is an isomorphism.

You may also consult the eigenvalues of matrices.

These are some of the fundamental matrix operations. There are other advanced operations as well, such as matrix inversion, determinant computation, and eigenvalue calculations, which are important in linear algebra and various applications of matrices.

Primitives of continuous functions

0

Primitives of continuous functions play a key role in the calculation of integrals. Therefore, in this article, we will use examples to introduce a simple technique for determining continuous function primitives.

We assume that the reader knows how to calculate integrals.

Generalities on primitives of continuous functions

Let $f:D\subset \mathbb{R}$ be a continuous function. The primitive of the function $f$ is a differentiable function $F:D\to\mathbb{R}$ such that $F'(x)=f(x)$ for any $x\in D$.

If $f$ admits a primitive $F$, then the function $G=F+c$ is a primitive of $f$ for any constant $c\in\mathbb{R}$. Two primitives of $f$ differ from a constant.

Let us now state the fundamental theorem of calculus: Let $f:D\to \mathbb{R}$ be a continuous function. The n for any $a\in D,$ the function defined by We also write \begin{align*}F(x)=\int^x_a f(t)dt,\quad x\in D,\end{align*} is differentiable on $D$ and $F’=f$.

Remark that the primitive a function of $C^1$ class. This means that the derivative function $F’$ is continuous.

Examples of how to find primitives 

To determine a primitive of a continuous function, you need to know some classical primitives as well as partial fraction decomposition rules.

  • Determine a primitive of the function \begin{align*} f(x)=\frac{2x^2-3x+4}{(x-1)^2}\textrm{ on }]1,+\infty[ .\end{align*} We look for real constants $a,b$ and $c$ such that \begin{align*} f(x)=a+\frac{b}{x-1}+\frac {c}{(x-1)^2}. \end{align*} By putting everything at the same denominator in the right-hand side, we find \begin{align*} f(x)=\frac{ax^2+x(-2a+b)+(a-b+c)}{(x-1)^2}. \end{align*} According to the first expression of $f,$ one obtains $a=2,;b=1$ and $c=3$, so that for any $x>1,$ \begin{align*} f(x)=2+\frac{1}{x-1}+\frac {3}{(x-1)^2}. \end{align*} The primitive of $f$ on $]1,+\infty[$ is then given by \begin{align*} F(x)=2x+\ln(x-1)-\frac {3}{x-1}+C,\qquad C\in \mathbb{R} \end{align*}
  • Determine the primitive of the function \begin{align*} f(x)=\frac{2x-1}{1+x^2}\textrm{ on }\mathbb{R}.\end{align*} We first observe that $f$ is defined on $\mathbb{R}$. On the other hand, we can write \begin{align*} f(x)&=\frac{2x}{1+x^2}-\frac{1}{1+x^2}\cr &=\frac{(1+x^2)’}{1+x^2}-\frac{1}{1+x^2}\cr &= \frac{d}{dx}\left(\ln(1+x^2)-\arctan(x)\right). \end{align*} Thus, the primitive of $f$ on $\mathbb{R}$ is given by \begin{align*} F(x)=\ln(1+x^2)-\arctan(x)+C,\qquad C\in\mathbb{R}. \end{align*}
  • Calculate the primitive of the following function\begin{align*} f(x)=x e^{\lambda x}\textrm{ on }\mathbb{R}.\end{align*} In this question, we will use integration by parts method. We then have for any $\lambda\neq 0,$ \begin{align*} \int^x_c t e^{\lambda t}dt&= \int^x_c \frac{t}{\lambda} \left(e^{\lambda t}\right)’dt\cr & =\left[\frac{t}{\lambda} e^{\lambda t}\right]^x_c-\int^x_c \frac{1}{\lambda}e^{\lambda t}dt. \end{align*} Hence the primitive of $f$ on $\mathbb{R}$ is given by \begin{align*} F(x)=\frac{1}{\lambda} xe^{\lambda x}-\frac{1}{\lambda^2} e^{\lambda x}+C,\qquad C\in\mathbb{R}. \end{align*}
  • Compute the primitive of the function\begin{align*} f(x)=\frac{1}{\sqrt{x^2+2x+5}}\textrm{ on }\mathbb{R}.\end{align*} Observe that $f$ is defined on all $\mathbb{R}$. Here we shall use the change of variables technique. We have \begin{align*} \int^x_c \frac{dt}{\sqrt{t^2+2t+5}}&=\int^x_c \frac{dt}{\sqrt{(t^2+2t+4)+1}}\cr &=\int^x_c \frac{dt}{\sqrt{(t+2)^2+1}}. \end{align*} We now put $u=t+2$ then $du=dt$ and \begin{align*} \int^x_c \frac{dt}{\sqrt{t^2+2t+5}}&=\int^{x+2}_{c+2} \frac{du}{\sqrt{u^2+1}}\cr &= 2\left[\ln(u+\sqrt{u^2+1})\right]^{x+2}_{c+2}\cr &= 2\ln(x+2+\sqrt{x^2+2x+5})+C. \end{align*} This implies that the primitive of $f$ is \begin{align*} F(x)=2\ln(x+2+\sqrt{x^2+2x+5})+C. \end{align*}